The AI Reader.
The Evidence

The diffusion paper behind today's image generators could not take a prompt

Ho, Jain, and Abbeel's 2020 paper reached a state-of-the-art score on the 32 by 32 CIFAR-10 set by learning to reverse a thousand steps of added noise.

Why this matters

Diffusion is the method behind the current wave of image generators, and one 2020 paper is where people trace it. Open that paper and the pictures inside come as a surprise: 32 by 32 thumbnails of horses and boats, generated from noise with no caption anywhere, each one costing a thousand sequential passes through the network. This lesson reads what Jonathan Ho, Ajay Jain, and Pieter Abbeel actually demonstrated and at what scale, then traces where the parts you would recognize, a typed prompt, a fast render, a photoreal face, entered the story. The paper founded a research line. The image generator you are picturing was built on top of it one capability at a time, by other papers, over the two years that followed.

Background optional reading
01
How an image generator works: how a model turns a field of noise into a picture by predicting the noise and removing it, the mechanism this paper introduced.
02
The GAN paper: the earlier approach to generating images, and the standard scores this paper drew even with and then passed.

The paper that made diffusion generate real pictures

“Denoising Diffusion Probabilistic Models” is a 2020 paper by Jonathan Ho, Ajay Jain, and Pieter Abbeel at UC Berkeley, presented at the NeurIPS 2020 conference.1 Its method is the one this course covers in the lesson on how an image generator works: a fixed forward process takes a real image and adds a little random noise to it, over and over, until the picture is pure static. A network then learns to run that chain backward, subtracting a little of the noise at each step until an image reappears.1

What made the paper matter is a claim its authors kept deliberately small. Diffusion models were not new in 2020. What no one had shown was that they could produce sharp, believable images, and this paper showed exactly that.1 Almost everything else people now attach to the phrase “diffusion model” came later, from other papers.

State of the art at 32 by 32

Start with what the paper put on the page. Its headline result is on CIFAR-10, a standard set of 32 by 32 color photographs of ten kinds of object.1 On that set the model reached a Fréchet Inception Distance of 3.17, a standard score for how close a batch of generated images is to real ones, where lower is better. It was the best any method had reported at the time.1 Both that figure and its companion Inception score of 9.46 are unconditional numbers: the model was given no label, no class, and no caption, and drew each image from noise alone.1

To read 3.17, set it beside the field. The paper's own table puts the diffusion model against the strongest GAN of the day, the adversarial approach an earlier lesson covers on its own, and against the score-based model it improved on most directly.1

CIFAR-10 results as the paper reports them; FID lower is better, Inception score higher is better.1
Model Inception FID
NCSN (score-based, 2019) 8.87 25.32
StyleGAN2 + ADA (GAN) 9.74 3.26
DDPM 9.46 3.17

On FID the diffusion model came first, 3.17 against the leading GAN's 3.26. On Inception score the GAN was still ahead, 9.74 against 9.46.1 The paper led on one of the two standard measures and trailed by a hair on the other. The score-based model from a year earlier sat at an FID of 25.32, an order of magnitude worse, and it is the closest predecessor the paper beat outright.1

That score came at a cost the paper is equally plain about. Producing one image means running the whole reverse chain, one network evaluation per step, and the paper set that chain to a thousand steps for every experiment.1 Sampling a batch of 128 images took about 300 seconds on the authors' 2020 hardware, near 2.3 seconds an image, and the network doing the work held about 35.7 million weights.1

1,000Network passes to make one image
~2.3 sPer image, on 2020 hardware
35.7MWeights in the CIFAR-10 model

The paper did reach past CIFAR-10, to 256 by 256 photographs of bedrooms, churches, cats, and faces, again unconditional.1 There the results were more mixed: an FID of 4.90 on bedrooms and 7.89 on churches, but 19.75 on cats. The abstract sums the whole 256 by 256 effort up as “sample quality similar to ProgressiveGAN,” a GAN from three years before.1 The ceiling was 256 pixels on a side, one class of image at a time.

Diffusion did not start here

The paper is often called the founding document of diffusion. It founded something narrower, and said so. The idea of building a generative model by slowly destroying the structure in data with a forward diffusion and then learning to reverse it was introduced five years earlier, in 2015, by Jascha Sohl-Dickstein and three coauthors.2 The specific technical footing the 2020 paper builds on came in 2019 from Yang Song and Stefano Ermon. They estimated the direction toward more realistic data at many noise levels and sampled by following it. Their model reached an Inception score of 8.87 on CIFAR-10, the mark this paper passed.3 Ho, Jain, and Abbeel credit both openly.

In the paper's words
To the best of our knowledge, there has been no demonstration that they are capable of generating high quality samples. We show that diffusion models actually are capable of generating high quality samples.1 Ho, Jain, and Abbeel, Section 1

So the contribution is exact. Not diffusion, but the first demonstration that diffusion generates high-quality images, together with the training recipe that got it there. That recipe is the part practitioners kept. Rather than have the network predict the cleaned-up image at each step, the paper rewrites the problem so the network predicts the noise that was added. It trains the network on a plain squared-error gap between the real noise and the predicted noise.1 Predicting the noise, on that simple objective, is the move nearly every diffusion model since has inherited.1 Calling the paper a founding document is fair for the modern wave of image generation. It is wrong for diffusion itself.

Where the prompt actually comes from

Now set the 2020 paper next to a modern text-to-image tool. There a typed sentence becomes a photoreal image at high resolution in a few seconds, on ordinary hardware. Every one of those parts, the sentence, the speed, the resolution, the hardware, comes from a paper published after this one.

The thousand-step wait was the first to go. Denoising Diffusion Implicit Models, from Jiaming Song, Chenlin Meng, and Stefano Ermon, appeared in October 2020, four months after the original paper. It keeps the trained network and changes only how the reverse chain is walked, producing comparable images 10 to 50 times faster.4 Fast sampling was not a later era. It arrived almost at once.

Running on ordinary hardware came from moving the work off the pixels. The 2020 model denoises the full image directly, which is what makes 256 pixels a strain. Latent Diffusion Models, from Robin Rombach and colleagues in 2021, run the same process in the compressed latent space of a pretrained autoencoder, a network trained to squeeze an image down to a small code and rebuild it. Working on that small code instead of the full picture cuts the computation enough to reach high resolution on modest hardware.5 That paper is the basis of Stable Diffusion, released to the public in August 2022 and able to run on a consumer graphics card with as little as 2.4 gigabytes of memory.6

The typed prompt came from the same Latent Diffusion paper, which added cross-attention layers, a mechanism that lets the denoising network read a text description while it works.5 Making the words actually bite took one more idea: classifier-free guidance, from Jonathan Ho and Tim Salimans in 2022. It trains a single model to run both with and without the caption, then pushes each sample toward the captioned version. The trade is some variety for images that follow the prompt more closely.7

Photoreal pictures from a caption were shown at scale by two 2022 systems. DALL-E 2, from OpenAI, first turns the caption into a compact numeric description of the image it should produce, then lets a diffusion model paint the picture from that description.8 Imagen, from Google, feeds the text through a large frozen language model and cascades diffusion models up to high resolution, reporting a COCO image-quality score of 7.27 on the same lower-is-better scale.9 These are the papers behind the products. The 2020 paper is behind these papers.

The takeaway

Here is the whole of it. In 2020, Ho, Jain, and Abbeel showed for the first time that a diffusion model could generate sharp images. They showed it on unconditional 32 by 32 CIFAR-10, where they set a state-of-the-art FID of 3.17. They showed it again on 256 by 256 bedrooms and faces, one class at a time, at a thousand slow steps an image. They did not invent diffusion, which is five years older, and they did not build anything you would call a product: no prompt, no speed, no photorealism, no laptop. Those came over the next two years, each in its own paper: fast sampling within months, then latent-space diffusion, text conditioning, guidance that makes a prompt bite, and the two big text-to-image systems. When someone calls this the paper that started AI image generation, the honest version is narrower. It is the paper that first made diffusion work, and everything you picture was built on that one result, by other people, afterward.

Go deeper optional reading
01
Lilian Weng, “What are Diffusion Models?”: the forward and reverse process worked through in full math, including the noise-prediction objective this paper introduced.
02
Jay Alammar, “The Illustrated Stable Diffusion”: how the later pieces, latent space and text conditioning, fit together in the tool most people actually use.

Sources

  1. arXiv · Ho, Jain & Abbeel, “Denoising Diffusion Probabilistic Models” (2020)
  2. arXiv · Sohl-Dickstein et al., “Deep Unsupervised Learning using Nonequilibrium Thermodynamics” (2015)
  3. arXiv · Song & Ermon, “Generative Modeling by Estimating Gradients of the Data Distribution” (2019)
  4. arXiv · Song, Meng & Ermon, “Denoising Diffusion Implicit Models” (2020)
  5. arXiv · Rombach et al., “High-Resolution Image Synthesis with Latent Diffusion Models” (2021)
  6. Wikipedia · Stable Diffusion (release date and hardware)
  7. arXiv · Ho & Salimans, “Classifier-Free Diffusion Guidance” (2022)
  8. arXiv · Ramesh et al., “Hierarchical Text-Conditional Image Generation with CLIP Latents” (2022)
  9. arXiv · Saharia et al., “Photorealistic Text-to-Image Diffusion Models with Deep Language Understanding” (2022)